diff --git a/img_styler/image_prompt/stable_diffusion.py b/img_styler/image_prompt/stable_diffusion.py
new file mode 100644
index 0000000..d8939d1
--- /dev/null
+++ b/img_styler/image_prompt/stable_diffusion.py
@@ -0,0 +1,27 @@
+import torch
+from torch import autocast
+from PIL import Image
+from io import BytesIO
+from diffusers import StableDiffusionImg2ImgPipeline
+
+
+def generate_image_with_prompt(input_img_path: str, prompt_txt: str = "Face portrait",
+                                                        output_path: str=None, save: bool=True):
+    torch.cuda.empty_cache()
+    device = "cuda"
+    # License: https://huggingface.co/spaces/CompVis/stable-diffusion-license
+    model_path = "./models/stable_diffusion_v1_4"
+    pipe = StableDiffusionImg2ImgPipeline.from_pretrained(model_path, revision="fp16",
+                                                                torch_dtype=torch.float16,)
+    pipe = pipe.to(device)
+
+    # Open image
+    image_input = Image.open(input_img_path).convert("RGB")
+    init_image = image_input.resize((512, 512))
+
+    with autocast(device):
+        images = pipe(prompt=prompt_txt, init_image=init_image, strength=0.5, guidance_scale=7.5)["sample"]
+
+    file_name = output_path + '/result.jpg'
+    images[0].save(file_name)
+    return file_name